--- tags: - text-to-image - stable-diffusion - lora - diffusers - template:sd-lora base_model: stabilityai/stable-diffusion-xl-base-1.0 license: cc-by-nc-nd-4.0 --- # ⚡ Flash Diffusion: FlashSDXL ⚡ Flash Diffusion is a diffusion distillation method proposed in [Flash Diffusion: Accelerating Any Conditional Diffusion Model for Few Steps Image Generation](http://arxiv.org/abs/2406.02347) *by Clément Chadebec, Onur Tasar, Eyal Benaroche, and Benjamin Aubin.* This model is a **108M** LoRA distilled version of [SDXL](https://huggingface.co./stabilityai/stable-diffusion-xl-base-1.0) model that is able to generate images in **4 steps**. The main purpose of this model is to reproduce the main results of the paper. See our [live demo](https://huggingface.co./spaces/jasperai/FlashPixart).

# How to use? The model can be used using the `DiffusionPipeline` from `diffusers` library directly. It can allow reducing the number of required sampling steps to **4 steps**. ```python from diffusers import DiffusionPipeline, LCMScheduler adapter_id = "jasperai/flash-sdxl" pipe = DiffusionPipeline.from_pretrained( "stabilityai/stable-diffusion-xl-base-1.0", use_safetensors=True, ) pipe.scheduler = LCMScheduler.from_pretrained( "stabilityai/stable-diffusion-xl-base-1.0", subfolder="scheduler", timestep_spacing="trailing", ) pipe.to("cuda") # Fuse and load LoRA weights pipe.load_lora_weights(adapter_id) pipe.fuse_lora() prompt = "A raccoon reading a book in a lush forest." image = pipe(prompt, num_inference_steps=4, guidance_scale=0).images[0] ```

# Training Details The model was trained for 20k iterations on 4 H100 GPUs (representing approximately a total of 176 GPU hours of training). Please refer to the [paper](http://arxiv.org/abs/2406.02347) for further parameters details. **Metrics on COCO 2014 validation (Table 3)** - FID-10k: 21.62 (4 NFE) - CLIP Score: 0.327 (4 NFE) ## License This model is released under the the Creative Commons BY-NC license.